What is Stable Diffusion XL 0.9 and How to Use It?

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In the rapidly evolving world of artificial intelligence, new models are constantly being developed and refined. One such model that has recently made waves in the AI community is the Stable Diffusion XL 0.9 (SDXL 0.9). This article delves into the details of SDXL 0.9, comparing it with other models in the Stable Diffusion series and the Midjourney V5 model. We’ll explore its unique features, advantages, and limitations, and provide a comprehensive guide on how to use it. Whether you’re an AI enthusiast, a researcher, or someone interested in the latest developments in AI, this article will provide you with valuable insights into the capabilities of SDXL 0.9.

Table of Contents

What is Stable Diffusion XL 0.9?

Stable Diffusion XL 0.9 (SDXL 0.9) is the latest development in Stability AI’s Stable Diffusion text-to-image suite of models. It’s a powerful AI tool capable of generating hyper-realistic creations for various applications, including films, television, music, instructional videos, and design and industrial use. It also offers functionalities beyond basic text prompting, such as image-to-image prompting, inpainting, and outpainting.

What's the Difference between Stable Diffusion XL 0.9 and Other Stable Diffusion Models?

The key difference between SDXL 0.9 and its predecessors lies in its significant increase in parameter count. SDXL 0.9 boasts a 3.5B parameter base model and a 6.6B parameter model ensemble pipeline. It runs on two CLIP models, including one of the largest OpenCLIP models trained to date, which enables it to create realistic imagery with greater depth and a higher resolution of 1024×1024. This is a significant improvement over the beta version, which runs on 3.1B parameters and uses just a single model. We can see two examples From SDXL Beta and SDXL 0.9: Image Source:stability.ai

Pros:

  • Improved Image Quality: SDXL 0.9 offers a significant improvement in image and composition detail over its predecessor.
  • Advanced Functionalities: Beyond basic text prompting, it offers functionalities like image-to-image prompting, inpainting, and outpainting.
  • Accessible: Despite its advanced capabilities, SDXL 0.9 can be run on a modern consumer GPU.

Cons:

  • Limited Availability: Currently, SDXL 0.9 is provided for research purposes only during a limited period.
  • Non-Commercial License: It is released under a non-commercial, research-only license.

SDXL 0.9 can be accessed via the ClipDrop platform, with an API coming soon. You’ll need a modern consumer GPU, a Windows 10 or 11, or Linux operating system, with 16GB RAM, and an Nvidia GeForce RTX 20 graphics card (or equivalent) with a minimum of 8GB of VRAM. The code to run it will soon be publicly available on Github.

Based on the information from the ClipDrop website, here’s a step-by-step guide on how to use Stable Diffusion XL 0.9:

  1. Access ClipDrop: Open your web browser and navigate to the ClipDrop website. If you haven’t already, you may need to create an account and sign in.
  2. Select Stable Diffusion XL 0.9: Once you’re logged in, navigate to the Stable Diffusion XL 0.9 page.
  3. Input Your Prompt: You will see a text box where you can input your prompt. This could be a description of the image you want the model to generate. Enter your prompt here.
  4. Click “Generate”: After entering your prompt, click the “Generate” button to start the image generation process.
  5. Wait for the Output: The model will take some time to generate the image. Wait for this process to complete.
  6. View Your Image: Once the image has been generated, it will be displayed on the screen. You can view the image to see if it matches your expectations.
  7. Save Your Image: If you’re satisfied with the generated image, you can save it. There should be an option to save or download the image. Click this button to save the image to your device.

Stable Diffusion XL 0.9 VS Midjourney V5

Midjourney V5 is another advanced AI model that produces more photographic generations than the default 5.1 model. This model produces images that closely match the prompt but may require longer prompts to achieve your desired aesthetic. It was the default model from Nov 2022–May 2023 and featured an entirely new codebase and brand-new AI architecture designed by Midjourney and trained on the new Midjourney AI supercluster. Model Version 4 had increased knowledge of creatures, places, and objects compared to previous models.

To further illustrate the differences and similarities between these two models, let’s take a look at the following comparison table:

Feature/Model Stable Diffusion XL 0.9 Midjourney V5
Parameter Count
3.5B base model, 6.6B ensemble pipeline
Not specified
Image Resolution(Default)
1024×1024
1115×625
Image Generation
Hyper-realistic images
More photographic generations
Additional Functionalities
Image-to-image prompting, inpainting, outpainting
Requires longer prompts for desired aesthetic
Release Date
June 2023
Default from Nov 2022–May 2023
Access

Via ClipDrop, API coming soon

License
Non-commercial, research-only
Not specified

FAQ

SDXL 0.9 can be used for various applications, including films, television, music, instructional videos, and design and industrial use. It also offers functionalities beyond basic text prompting, such as image-to-image prompting, inpainting, and outpainting.

Currently, SDXL 0.9 is provided for research purposes only during a limited period. It is released under a non-commercial, research-only license.

The key difference lies in its significant increase in parameter count and the use of two CLIP models, which enables it to create realistic imagery with greater depth and a higher resolution of 1024×1024.

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